We present the Recurrent Interface Network (RIN), a neural net architecture that allocates computation adaptively to the input according to the distribution of information, allowing it to scale to iterative generation of high-dimensional data. Hidden units of RINs are partitioned into the interface, which is locally connected to inputs, and latents, which are decoupled from inputs and can exchange information globally. The RIN block selectively reads from the interface into latents for high-capacity processing, with incremental updates written back to the interface. Stacking multiple blocks enables effective routing across local and global levels. While routing adds overhead, the cost can be amortized in recurrent computation settings where inputs change gradually while more global context persists, such as iterative generation using diffusion models. To this end, we propose a latent self-conditioning technique that "warm-starts" the latents at each iteration of the generation process. When applied to diffusion models operating directly on pixels, RINs yield state-of-the-art image and video generation without cascades or guidance, while being domain-agnostic and up to 10$\times$ more efficient compared to specialized 2D and 3D U-Nets.

We explore a new class of diffusion models based on the transformer architecture. We train latent diffusion models of images, replacing the commonly-used U-Net backbone with a transformer that operates on latent patches. We analyze the scalability of our Diffusion Transformers (DiTs) through the lens of forward pass complexity as measured by Gflops. We find that DiTs with higher Gflops -- through increased transformer depth/width or increased number of input tokens -- consistently have lower FID. In addition to possessing good scalability properties, our largest DiT-XL/2 models outperform all prior diffusion models on the class-conditional ImageNet 512x512 and 256x256 benchmarks, achieving a state-of-the-art FID of 2.27 on the latter.

我们提出了位扩散：一种简单而通用的方法，用于通过连续扩散模型生成离散数据。我们方法背后的主要思想是首先将离散数据表示为二进制位，然后训练连续扩散模型，以将这些位模拟为实数，我们称为模拟位。为了生成样品，模型首先生成模拟位，然后将其阈值阈值以获得表示离散变量的位。我们进一步提出了两种简单的技术，即自我调节和不对称的时间间隔，从而导致样本质量的显着改善。尽管它很简单，但提出的方法可以在离散图像生成和图像字幕任务中实现强大的性能。对于离散图像生成，我们在CIFAR-10（具有3K离散的8位代币）和Imagenet-64x64（具有12K离散的8位代币）上都显着改善了先前的最新技术，超过了最好的自动回归性自动回应。样品质量（通过FID衡量）和效率的模型。对于MS-Coco数据集上的图像字幕，与自回归模型相比，我们的方法可实现竞争成果。

Score-based diffusion models have captured widespread attention and funded fast progress of recent vision generative tasks. In this paper, we focus on diffusion model backbone which has been much neglected before. We systematically explore vision Transformers as diffusion learners for various generative tasks. With our improvements the performance of vanilla ViT-based backbone (IU-ViT) is boosted to be on par with traditional U-Net-based methods. We further provide a hypothesis on the implication of disentangling the generative backbone as an encoder-decoder structure and show proof-of-concept experiments verifying the effectiveness of a stronger encoder for generative tasks with ASymmetriC ENcoder Decoder (ASCEND). Our improvements achieve competitive results on CIFAR-10, CelebA, LSUN, CUB Bird and large-resolution text-to-image tasks. To the best of our knowledge, we are the first to successfully train a single diffusion model on text-to-image task beyond 64x64 resolution. We hope this will motivate people to rethink the modeling choices and the training pipelines for diffusion-based generative models.

扩散模型是一类强大的生成模型类别，可以迭代地贬低样品生成数据。尽管许多作品都集中在此抽样过程中的迭代次数上，但很少有人专注于每次迭代的成本。我们发现，添加简单的VIT风格的修补转换可以大大减少扩散模型的采样时间和内存使用情况。我们通过对扩散模型目标的分析以及在LSUN教堂，Imagenet 256和FFHQ1024上进行的经验实验来证明我们的方法是合理的。我们在Tensorflow和Pytorch中提供了实现。

We show that diffusion models can achieve image sample quality superior to the current state-of-the-art generative models. We achieve this on unconditional image synthesis by finding a better architecture through a series of ablations. For conditional image synthesis, we further improve sample quality with classifier guidance: a simple, compute-efficient method for trading off diversity for fidelity using gradients from a classifier. We achieve an FID of 2.97 on ImageNet 128×128, 4.59 on ImageNet 256×256, and 7.72 on ImageNet 512×512, and we match BigGAN-deep even with as few as 25 forward passes per sample, all while maintaining better coverage of the distribution. Finally, we find that classifier guidance combines well with upsampling diffusion models, further improving FID to 3.94 on ImageNet 256×256 and 3.85 on ImageNet 512×512. We release our code at https://github.com/openai/guided-diffusion.

Panoptic segmentation assigns semantic and instance ID labels to every pixel of an image. As permutations of instance IDs are also valid solutions, the task requires learning of high-dimensional one-to-many mapping. As a result, state-of-the-art approaches use customized architectures and task-specific loss functions. We formulate panoptic segmentation as a discrete data generation problem, without relying on inductive bias of the task. A diffusion model based on analog bits is used to model panoptic masks, with a simple, generic architecture and loss function. By simply adding past predictions as a conditioning signal, our method is capable of modeling video (in a streaming setting) and thereby learns to track object instances automatically. With extensive experiments, we demonstrate that our generalist approach can perform competitively to state-of-the-art specialist methods in similar settings.

DeNoising扩散模型代表了计算机视觉中最新的主题，在生成建模领域表现出了显着的结果。扩散模型是一个基于两个阶段的深层生成模型，一个正向扩散阶段和反向扩散阶段。在正向扩散阶段，通过添加高斯噪声，输入数据在几个步骤中逐渐受到干扰。在反向阶段，模型的任务是通过学习逐步逆转扩散过程来恢复原始输入数据。尽管已知的计算负担，即由于采样过程中涉及的步骤数量，扩散模型对生成样品的质量和多样性得到了广泛赞赏。在这项调查中，我们对视觉中应用的denoising扩散模型的文章进行了全面综述，包括该领域的理论和实际贡献。首先，我们识别并介绍了三个通用扩散建模框架，这些框架基于扩散概率模型，噪声调节得分网络和随机微分方程。我们进一步讨论了扩散模型与其他深层生成模型之间的关系，包括变异自动编码器，生成对抗网络，基于能量的模型，自回归模型和正常流量。然后，我们介绍了计算机视觉中应用的扩散模型的多角度分类。最后，我们说明了扩散模型的当前局限性，并设想了一些有趣的未来研究方向。

Denoising diffusion probabilistic models (DDPM) are a class of generative models which have recently been shown to produce excellent samples. We show that with a few simple modifications, DDPMs can also achieve competitive loglikelihoods while maintaining high sample quality. Additionally, we find that learning variances of the reverse diffusion process allows sampling with an order of magnitude fewer forward passes with a negligible difference in sample quality, which is important for the practical deployment of these models. We additionally use precision and recall to compare how well DDPMs and GANs cover the target distribution. Finally, we show that the sample quality and likelihood of these models scale smoothly with model capacity and training compute, making them easily scalable. We release our code at https://github.com/ openai/improved-diffusion.

生成建模研究的持续趋势是将样本分辨率推高更高，同时减少培训和采样的计算要求。我们的目标是通过技术的组合进一步推动这一趋势 - 每个组件代表当前效率在各自领域的顶峰。其中包括载体定量的GAN（VQ-GAN），该模型具有高水平的损耗 - 但感知上微不足道的压缩模型；沙漏变形金刚，一个高度可扩展的自我注意力模型；和逐步未胶片的denoising自动编码器（Sundae），一种非自动化（NAR）文本生成模型。出乎意料的是，当应用于多维数据时，我们的方法突出了沙漏变压器的原始公式中的弱点。鉴于此，我们建议对重采样机制进行修改，该机制适用于将分层变压器应用于多维数据的任何任务。此外，我们证明了圣代表到长序列长度的可伸缩性 - 比先前的工作长四倍。我们提出的框架秤达到高分辨率（$ 1024 \ times 1024 $），并迅速火车（2-4天）。至关重要的是，训练有素的模型在消费级GPU（GTX 1080TI）上大约2秒内生产多样化和现实的百像样品。通常，该框架是灵活的：支持任意数量的采样步骤，示例自动插入，自我纠正功能，有条件的生成和NAR公式，以允许任意介绍掩护。我们在FFHQ256上获得10.56的FID得分 - 仅在100个采样步骤中以不到一半的采样步骤接近原始VQ -GAN，而FFHQ1024的FFHQ1024和21.85。

利用深度学习的最新进展，文本到图像生成模型目前具有吸引公众关注的优点。其中两个模型Dall-E 2和Imagen已经证明，可以从图像的简单文本描述中生成高度逼真的图像。基于一种称为扩散模型的新型图像生成方法，文本对图像模型可以生产许多不同类型的高分辨率图像，其中人类想象力是唯一的极限。但是，这些模型需要大量的计算资源来训练，并处理从互联网收集的大量数据集。此外，代码库和模型均未发布。因此，它可以防止AI社区尝试这些尖端模型，从而使其结果复制变得复杂，即使不是不可能。在本文中，我们的目标是首先回顾这些模型使用的不同方法和技术，然后提出我们自己的文本模型模型实施。高度基于DALL-E 2，我们引入了一些轻微的修改，以应对所引起的高计算成本。因此，我们有机会进行实验，以了解这些模型的能力，尤其是在低资源制度中。特别是，我们提供了比Dall-e 2的作者（包括消融研究）更深入的分析。此外，扩散模型使用所谓的指导方法来帮助生成过程。我们引入了一种新的指导方法，该方法可以与其他指导方法一起使用，以提高图像质量。最后，我们的模型产生的图像质量相当好，而不必维持最先进的文本对图像模型的重大培训成本。

Large-scale diffusion-based generative models have led to breakthroughs in text-conditioned high-resolution image synthesis. Starting from random noise, such text-to-image diffusion models gradually synthesize images in an iterative fashion while conditioning on text prompts. We find that their synthesis behavior qualitatively changes throughout this process: Early in sampling, generation strongly relies on the text prompt to generate text-aligned content, while later, the text conditioning is almost entirely ignored. This suggests that sharing model parameters throughout the entire generation process may not be ideal. Therefore, in contrast to existing works, we propose to train an ensemble of text-to-image diffusion models specialized for different synthesis stages. To maintain training efficiency, we initially train a single model, which is then split into specialized models that are trained for the specific stages of the iterative generation process. Our ensemble of diffusion models, called eDiff-I, results in improved text alignment while maintaining the same inference computation cost and preserving high visual quality, outperforming previous large-scale text-to-image diffusion models on the standard benchmark. In addition, we train our model to exploit a variety of embeddings for conditioning, including the T5 text, CLIP text, and CLIP image embeddings. We show that these different embeddings lead to different behaviors. Notably, the CLIP image embedding allows an intuitive way of transferring the style of a reference image to the target text-to-image output. Lastly, we show a technique that enables eDiff-I's "paint-with-words" capability. A user can select the word in the input text and paint it in a canvas to control the output, which is very handy for crafting the desired image in mind. The project page is available at https://deepimagination.cc/eDiff-I/

现实世界中的数据是高维的：即使在压缩后，书籍，图像或音乐表演也很容易包含数十万个元素。但是，最常用的自回归模型，变压器非常昂贵，以缩放捕获这种远程结构所需的输入和层数。我们开发了感知者AR，这是一种自回归的模态 - 不合骨架构，它使用交叉注意力将远程输入映射到少数潜在的潜在，同时还可以维护端到端的因果关系掩盖。感知器AR可以直接进行十万个令牌，从而实现了实用的长篇小写密度估计，而无需手工制作的稀疏模式或记忆机制。当对图像或音乐进行培训时，感知器AR会生成具有清晰长期连贯性和结构的输出。我们的架构还获得了长期基准测试的最新可能性，包括64 x 64个Imagenet图像和PG-19书籍。

最近，Rissanen等人（2022年）提出了一种基于热量耗散或模糊的生成建模的新型扩散过程，作为各向同性高斯扩散的替代方法。在这里，我们表明，可以通过与非各向异性噪声的高斯扩散过程来等效地定义模糊。在建立这一联系时，我们弥合了反热量耗散和降解扩散之间的缝隙，并阐明了由于这种建模选择而导致的感应偏置。最后，我们提出了一类普遍的扩散模型，该模型既可以提供标准的高斯denoisis扩散和逆热散热，我们称之为模糊的扩散模型。

Can continuous diffusion models bring the same performance breakthrough on natural language they did for image generation? To circumvent the discrete nature of text data, we can simply project tokens in a continuous space of embeddings, as is standard in language modeling. We propose Self-conditioned Embedding Diffusion, a continuous diffusion mechanism that operates on token embeddings and allows to learn flexible and scalable diffusion models for both conditional and unconditional text generation. Through qualitative and quantitative evaluation, we show that our text diffusion models generate samples comparable with those produced by standard autoregressive language models - while being in theory more efficient on accelerator hardware at inference time. Our work paves the way for scaling up diffusion models for text, similarly to autoregressive models, and for improving performance with recent refinements to continuous diffusion.

Biological systems perceive the world by simultaneously processing high-dimensional inputs from modalities as diverse as vision, audition, touch, proprioception, etc. The perception models used in deep learning on the other hand are designed for individual modalities, often relying on domainspecific assumptions such as the local grid structures exploited by virtually all existing vision models. These priors introduce helpful inductive biases, but also lock models to individual modalities. In this paper we introduce the Perceiver -a model that builds upon Transformers and hence makes few architectural assumptions about the relationship between its inputs, but that also scales to hundreds of thousands of inputs, like ConvNets. The model leverages an asymmetric attention mechanism to iteratively distill inputs into a tight latent bottleneck, allowing it to scale to handle very large inputs. We show that this architecture is competitive with or outperforms strong, specialized models on classification tasks across various modalities: images, point clouds, audio, video, and video+audio. The Perceiver obtains performance comparable to ResNet-50 and ViT on ImageNet without 2D convolutions by directly attending to 50,000 pixels. It is also competitive in all modalities in AudioSet.

我们使用条件扩散模型介绍调色板，这是一种简单而一般的框架，可用于图像到图像到图像转换。在四个具有挑战性的图像到图像转换任务（着色，染色，un折叠和JPEG减压），调色板优于强大的GaN和回归基线，并建立了新的最新状态。这是在没有特定于任务特定的超参数调整，架构定制或任何辅助损耗的情况下实现的，展示了理想的一般性和灵活性。我们揭示了使用$ l_2 $与vs. $ l_1 $损失在样本多样性上的越来越多的影响，并通过经验架构研究表明自我关注的重要性。重要的是，我们倡导基于想象项目的统一评估协议，并报告包括预先训练的Reset-50的FID，成立得分，分类准确度的多个样本质量评分，以及针对各种基线的参考图像的感知距离。我们预计这一标准化评估协议在推进图像到图像翻译研究方面发挥着关键作用。最后，我们表明，在3个任务（着色，染色，JPEG减压）上培训的单个通用调色板模型也表现或优于特定于任务专家的专家对应物。

Diffusion models have quickly become the go-to paradigm for generative modelling of perceptual signals (such as images and sound) through iterative refinement. Their success hinges on the fact that the underlying physical phenomena are continuous. For inherently discrete and categorical data such as language, various diffusion-inspired alternatives have been proposed. However, the continuous nature of diffusion models conveys many benefits, and in this work we endeavour to preserve it. We propose CDCD, a framework for modelling categorical data with diffusion models that are continuous both in time and input space. We demonstrate its efficacy on several language modelling tasks.

Diffusion models have achieved great success in modeling continuous data modalities such as images, audio, and video, but have seen limited use in discrete domains such as language. Recent attempts to adapt diffusion to language have presented diffusion as an alternative to autoregressive language generation. We instead view diffusion as a complementary method that can augment the generative capabilities of existing pre-trained language models. We demonstrate that continuous diffusion models can be learned in the latent space of a pre-trained encoder-decoder model, enabling us to sample continuous latent representations that can be decoded into natural language with the pre-trained decoder. We show that our latent diffusion models are more effective at sampling novel text from data distributions than a strong autoregressive baseline and also enable controllable generation.

通过将图像形成过程分解成逐个申请的去噪自身额，扩散模型（DMS）实现了最先进的合成导致图像数据和超越。另外，它们的配方允许引导机构来控制图像生成过程而不会再刷新。然而，由于这些模型通常在像素空间中直接操作，因此强大的DMS的优化通常消耗数百个GPU天，并且由于顺序评估，推理是昂贵的。为了在保留其质量和灵活性的同时启用有限计算资源的DM培训，我们将它们应用于强大的佩带自动化器的潜在空间。与以前的工作相比，这种代表上的培训扩散模型允许第一次达到复杂性降低和细节保存之间的近乎最佳点，极大地提高了视觉保真度。通过将跨关注层引入模型架构中，我们将扩散模型转化为强大而柔性的发电机，以进行诸如文本或边界盒和高分辨率合成的通用调节输入，以卷积方式变得可以实现。我们的潜在扩散模型（LDMS）实现了一种新的技术状态，可在各种任务中进行图像修复和高竞争性能，包括无条件图像生成，语义场景合成和超级分辨率，同时与基于像素的DMS相比显着降低计算要求。代码可在https://github.com/compvis/lattent-diffusion获得。