Recent 3D generative models have achieved remarkable performance in synthesizing high resolution photorealistic images with view consistency and detailed 3D shapes, but training them for diverse domains is challenging since it requires massive training images and their camera distribution information. Text-guided domain adaptation methods have shown impressive performance on converting the 2D generative model on one domain into the models on other domains with different styles by leveraging the CLIP (Contrastive Language-Image Pre-training), rather than collecting massive datasets for those domains. However, one drawback of them is that the sample diversity in the original generative model is not well-preserved in the domain-adapted generative models due to the deterministic nature of the CLIP text encoder. Text-guided domain adaptation will be even more challenging for 3D generative models not only because of catastrophic diversity loss, but also because of inferior text-image correspondence and poor image quality. Here we propose DATID-3D, a domain adaptation method tailored for 3D generative models using text-to-image diffusion models that can synthesize diverse images per text prompt without collecting additional images and camera information for the target domain. Unlike 3D extensions of prior text-guided domain adaptation methods, our novel pipeline was able to fine-tune the state-of-the-art 3D generator of the source domain to synthesize high resolution, multi-view consistent images in text-guided targeted domains without additional data, outperforming the existing text-guided domain adaptation methods in diversity and text-image correspondence. Furthermore, we propose and demonstrate diverse 3D image manipulations such as one-shot instance-selected adaptation and single-view manipulated 3D reconstruction to fully enjoy diversity in text.

Can a text-to-image diffusion model be used as a training objective for adapting a GAN generator to another domain? In this paper, we show that the classifier-free guidance can be leveraged as a critic and enable generators to distill knowledge from large-scale text-to-image diffusion models. Generators can be efficiently shifted into new domains indicated by text prompts without access to groundtruth samples from target domains. We demonstrate the effectiveness and controllability of our method through extensive experiments. Although not trained to minimize CLIP loss, our model achieves equally high CLIP scores and significantly lower FID than prior work on short prompts, and outperforms the baseline qualitatively and quantitatively on long and complicated prompts. To our best knowledge, the proposed method is the first attempt at incorporating large-scale pre-trained diffusion models and distillation sampling for text-driven image generator domain adaptation and gives a quality previously beyond possible. Moreover, we extend our work to 3D-aware style-based generators and DreamBooth guidance.

Recent CLIP-guided 3D optimization methods, e.g., DreamFields and PureCLIPNeRF achieve great success in zero-shot text-guided 3D synthesis. However, due to the scratch training and random initialization without any prior knowledge, these methods usually fail to generate accurate and faithful 3D structures that conform to the corresponding text. In this paper, we make the first attempt to introduce the explicit 3D shape prior to CLIP-guided 3D optimization methods. Specifically, we first generate a high-quality 3D shape from input texts in the text-to-shape stage as the 3D shape prior. We then utilize it as the initialization of a neural radiance field and then optimize it with the full prompt. For the text-to-shape generation, we present a simple yet effective approach that directly bridges the text and image modalities with a powerful text-to-image diffusion model. To narrow the style domain gap between images synthesized by the text-to-image model and shape renderings used to train the image-to-shape generator, we further propose to jointly optimize a learnable text prompt and fine-tune the text-to-image diffusion model for rendering-style image generation. Our method, namely, Dream3D, is capable of generating imaginative 3D content with better visual quality and shape accuracy than state-of-the-art methods.

最近，GaN反演方法与对比语言 - 图像预先绘制（CLIP）相结合，可以通过文本提示引导零拍摄图像操作。然而，由于GaN反转能力有限，它们对不同实物的不同实物的应用仍然困难。具体地，这些方法通常在与训练数据相比，改变对象标识或产生不需要的图像伪影的比较与新颖姿势，视图和高度可变内容重建具有新颖姿势，视图和高度可变内容的困难。为了减轻这些问题并实现真实图像的忠实操纵，我们提出了一种新的方法，Dumbused Clip，其使用扩散模型执行文本驱动的图像操纵。基于近期扩散模型的完整反转能力和高质量的图像生成功率，即使在看不见的域之间也成功地执行零拍摄图像操作。此外，我们提出了一种新颖的噪声组合方法，允许简单的多属性操作。与现有基线相比，广泛的实验和人类评估确认了我们的方法的稳健和卓越的操纵性能。

数字艺术合成在多媒体社区中受到越来越多的关注，因为有效地与公众参与了艺术。当前的数字艺术合成方法通常使用单模式输入作为指导，从而限制了模型的表现力和生成结果的多样性。为了解决这个问题，我们提出了多模式引导的艺术品扩散（MGAD）模型，该模型是一种基于扩散的数字艺术品生成方法，它利用多模式提示作为控制无分类器扩散模型的指导。此外，对比度语言图像预处理（剪辑）模型用于统一文本和图像模式。关于生成的数字艺术绘画质量和数量的广泛实验结果证实了扩散模型和多模式指导的组合有效性。代码可从https://github.com/haha-lisa/mgad-multimodal-guided-artwork-diffusion获得。

可以训练生成模型，以从特定域中生成图像，仅由文本提示引导，而不看到任何图像？换句话说：可以将图像生成器“盲目地训练”吗？利用大规模对比语言图像预训练（CLIP）模型的语义力量，我们提出了一种文本驱动方法，允许将生成模型转移到新域，而无需收集单个图像。我们展示通过自然语言提示和几分钟的培训，我们的方法可以通过各种风格和形状的多种域调整发电机。值得注意的是，许多这些修改难以与现有方法达到困难或完全不可能。我们在广泛的域中进行了广泛的实验和比较。这些证明了我们方法的有效性，并表明我们的移动模型保持了对下游任务吸引的生成模型的潜在空间属性。

可控图像合成模型允许根据文本指令或来自示例图像的指导创建不同的图像。最近，已经显示出去噪扩散概率模型比现有方法产生更现实的图像，并且已在无条件和类条件设置中成功展示。我们探索细粒度，连续控制该模型类，并引入了一种新颖的统一框架，用于语义扩散指导，允许语言或图像指导，或两者。使用图像文本或图像匹配分数的梯度将指导注入预训练的无条件扩散模型中。我们探讨基于剪辑的文本指导，以及以统一形式的基于内容和类型的图像指导。我们的文本引导综合方法可以应用于没有相关文本注释的数据集。我们对FFHQ和LSUN数据集进行实验，并显示出细粒度的文本引导图像合成的结果，与样式或内容示例图像相关的图像的合成，以及具有文本和图像引导的示例。

随着信息中的各种方式存在于现实世界中的各种方式，多式联信息之间的有效互动和融合在计算机视觉和深度学习研究中的多模式数据的创造和感知中起着关键作用。通过卓越的功率，在多式联运信息中建模互动，多式联运图像合成和编辑近年来已成为一个热门研究主题。与传统的视觉指导不同，提供明确的线索，多式联路指南在图像合成和编辑方面提供直观和灵活的手段。另一方面，该领域也面临着具有固有的模态差距的特征的几个挑战，高分辨率图像的合成，忠实的评估度量等。在本调查中，我们全面地阐述了最近多式联运图像综合的进展根据数据模型和模型架构编辑和制定分类。我们从图像合成和编辑中的不同类型的引导方式开始介绍。然后，我们描述了多模式图像综合和编辑方法，其具有详细的框架，包括生成的对抗网络（GAN），GaN反转，变压器和其他方法，例如NERF和扩散模型。其次是在多模式图像合成和编辑中广泛采用的基准数据集和相应的评估度量的综合描述，以及分析各个优点和限制的不同合成方法的详细比较。最后，我们为目前的研究挑战和未来的研究方向提供了深入了解。与本调查相关的项目可在HTTPS://github.com/fnzhan/mise上获得

大型文本对图像模型在AI的演变中取得了显着的飞跃，从而使图像从给定的文本提示中实现了高质量和多样化的图像合成。但是，这些模型缺乏在给定的参考集中模仿受试者的外观，并在不同情况下合成它们的新颖性。在这项工作中，我们提出了一种新的方法，用于“个性化”文本图像扩散模型（将它们专门针对用户的需求）。仅作为一个主题的几张图像给出，我们将验证的文本对图像模型（图像，尽管我们的方法不限于特定模型），以便它学会了将唯一标识符与该特定主题结合。一旦将受试者嵌入模型的输出域中，就可以使用唯一标识符来合成主题的完全新颖的光真逼真的图像在不同场景中的上下文化。通过利用具有新的自动构基特异性的先前保存损失的语义先验嵌入到模型中，我们的技术可以在参考图像中未出现的不同场景，姿势，视图和照明条件中合成主题。我们将技术应用于几个以前无用的任务，包括主题重新定义，文本指导的视图合成，外观修改和艺术渲染（所有这些都保留了主题的关键特征）。项目页面：https：//dreambooth.github.io/

Recent works on diffusion models have demonstrated a strong capability for conditioning image generation, e.g., text-guided image synthesis. Such success inspires many efforts trying to use large-scale pre-trained diffusion models for tackling a challenging problem--real image editing. Works conducted in this area learn a unique textual token corresponding to several images containing the same object. However, under many circumstances, only one image is available, such as the painting of the Girl with a Pearl Earring. Using existing works on fine-tuning the pre-trained diffusion models with a single image causes severe overfitting issues. The information leakage from the pre-trained diffusion models makes editing can not keep the same content as the given image while creating new features depicted by the language guidance. This work aims to address the problem of single-image editing. We propose a novel model-based guidance built upon the classifier-free guidance so that the knowledge from the model trained on a single image can be distilled into the pre-trained diffusion model, enabling content creation even with one given image. Additionally, we propose a patch-based fine-tuning that can effectively help the model generate images of arbitrary resolution. We provide extensive experiments to validate the design choices of our approach and show promising editing capabilities, including changing style, content addition, and object manipulation. The code is available for research purposes at https://github.com/zhang-zx/SINE.git .

文本对图像模型提供了前所未有的自由，可以通过自然语言指导创作。然而，尚不清楚如何行使这种自由以生成特定独特概念，修改其外观或以新角色和新颖场景构成它们的图像。换句话说，我们问：我们如何使用语言指导的模型将猫变成绘画，或者想象基于我们喜欢的玩具的新产品？在这里，我们提出了一种简单的方法，可以允许这种创造性自由。我们仅使用3-5个用户提供的概念（例如对象或样式）的图像，我们学会通过在冷冻文本到图像模型的嵌入空间中通过新的“单词”表示它。这些“单词”可以组成自然语言句子，以直观的方式指导个性化的创作。值得注意的是，我们发现有证据表明单词嵌入足以捕获独特而多样的概念。我们将我们的方法比较了各种基线，并证明它可以更忠实地描绘出一系列应用程序和任务的概念。我们的代码，数据和新单词将在以下网址提供：https：//textual-inversion.github.io

This paper presents a 3D generative model that uses diffusion models to automatically generate 3D digital avatars represented as neural radiance fields. A significant challenge in generating such avatars is that the memory and processing costs in 3D are prohibitive for producing the rich details required for high-quality avatars. To tackle this problem we propose the roll-out diffusion network (Rodin), which represents a neural radiance field as multiple 2D feature maps and rolls out these maps into a single 2D feature plane within which we perform 3D-aware diffusion. The Rodin model brings the much-needed computational efficiency while preserving the integrity of diffusion in 3D by using 3D-aware convolution that attends to projected features in the 2D feature plane according to their original relationship in 3D. We also use latent conditioning to orchestrate the feature generation for global coherence, leading to high-fidelity avatars and enabling their semantic editing based on text prompts. Finally, we use hierarchical synthesis to further enhance details. The 3D avatars generated by our model compare favorably with those produced by existing generative techniques. We can generate highly detailed avatars with realistic hairstyles and facial hair like beards. We also demonstrate 3D avatar generation from image or text as well as text-guided editability.

使用单视图2D照片仅集合，无监督的高质量多视图 - 一致的图像和3D形状一直是一个长期存在的挑战。现有的3D GAN是计算密集型的，也是没有3D-一致的近似;前者限制了所生成的图像的质量和分辨率，并且后者对多视图一致性和形状质量产生不利影响。在这项工作中，我们提高了3D GAN的计算效率和图像质量，而无需依赖这些近似。为此目的，我们介绍了一种表现力的混合明确隐式网络架构，与其他设计选择一起，不仅可以实时合成高分辨率多视图一致图像，而且还产生高质量的3D几何形状。通过解耦特征生成和神经渲染，我们的框架能够利用最先进的2D CNN生成器，例如Stylega2，并继承它们的效率和表现力。在其他实验中，我们展示了与FFHQ和AFHQ猫的最先进的3D感知合成。

我们将神经渲染与多模态图像和文本表示相结合，以仅从自然语言描述中综合不同的3D对象。我们的方法，梦场，可以产生多种物体的几何和颜色而无需3D监控。由于不同，标题3D数据的稀缺性，先前的方法仅生成来自少数类别的对象，例如ShapEnet。相反，我们指导生成与从Web的标题图像的大型数据集预先培训的图像文本模型。我们的方法优化了许多相机视图的神经辐射场，使得根据预先训练的剪辑模型，渲染图像非常高度地使用目标字幕。为了提高保真度和视觉质量，我们引入简单的几何前瞻，包括突出透射率正则化，场景界限和新的MLP架构。在实验中，梦场从各种自然语言标题中产生现实，多视图一致的物体几何和颜色。

Virtual reality and augmented reality (XR) bring increasing demand for 3D content. However, creating high-quality 3D content requires tedious work that a human expert must do. In this work, we study the challenging task of lifting a single image to a 3D object and, for the first time, demonstrate the ability to generate a plausible 3D object with 360{\deg} views that correspond well with the given reference image. By conditioning on the reference image, our model can fulfill the everlasting curiosity for synthesizing novel views of objects from images. Our technique sheds light on a promising direction of easing the workflows for 3D artists and XR designers. We propose a novel framework, dubbed NeuralLift-360, that utilizes a depth-aware neural radiance representation (NeRF) and learns to craft the scene guided by denoising diffusion models. By introducing a ranking loss, our NeuralLift-360 can be guided with rough depth estimation in the wild. We also adopt a CLIP-guided sampling strategy for the diffusion prior to provide coherent guidance. Extensive experiments demonstrate that our NeuralLift-360 significantly outperforms existing state-of-the-art baselines. Project page: https://vita-group.github.io/NeuralLift-360/

Text-guided diffusion models have shown superior performance in image/video generation and editing. While few explorations have been performed in 3D scenarios. In this paper, we discuss three fundamental and interesting problems on this topic. First, we equip text-guided diffusion models to achieve $\textbf{3D-consistent generation}$. Specifically, we integrate a NeRF-like neural field to generate low-resolution coarse results for a given camera view. Such results can provide 3D priors as condition information for the following diffusion process. During denoising diffusion, we further enhance the 3D consistency by modeling cross-view correspondences with a novel two-stream (corresponding to two different views) asynchronous diffusion process. Second, we study $\textbf{3D local editing}$ and propose a two-step solution that can generate 360$^{\circ}$ manipulated results by editing an object from a single view. Step 1, we propose to perform 2D local editing by blending the predicted noises. Step 2, we conduct a noise-to-text inversion process that maps 2D blended noises into the view-independent text embedding space. Once the corresponding text embedding is obtained, 360$^{\circ}$ images can be generated. Last but not least, we extend our model to perform \textbf{one-shot novel view synthesis} by fine-tuning on a single image, firstly showing the potential of leveraging text guidance for novel view synthesis. Extensive experiments and various applications show the prowess of our 3DDesigner. The project page is available at https://3ddesigner-diffusion.github.io/.

我们提出了快速的文本2stylegan，这是一种自然语言界面，可适应预先训练的甘体，以实现文本引导的人脸合成。利用对比性语言图像预训练（剪辑）的最新进展，在培训过程中不需要文本数据。Fast Text2Stylegan被配制为条件变异自动编码器（CVAE），可在测试时为生成的图像提供额外的控制和多样性。我们的模型在遇到新的文本提示时不需要重新训练或微调剂或剪辑。与先前的工作相反，我们不依赖于测试时间的优化，这使我们的方法数量级比先前的工作快。从经验上讲，在FFHQ数据集上，我们的方法提供了与先前的工作相比，自然语言描述中具有不同详细程度的自然语言描述中的图像。

最近已被证明扩散模型产生高质量的合成图像，尤其是与指导技术配对，以促进忠诚的多样性。我们探索文本条件图像综合问题的扩散模型，并比较了两种不同的指导策略：剪辑指导和自由分类指导。我们发现后者是人类评估者的优选，用于光敏和标题相似度，并且通常产生光素质拟种样品。使用自由分类指导的35亿参数文本条件扩散模型的样本由人类评估者对来自Dall-E的人的人们青睐，即使后者使用昂贵的剪辑重新划分。此外，我们发现我们的模型可以进行微调，以执行图像修复，从而实现强大的文本驱动的图像编辑。我们在过滤的数据集中培训较小的模型，并在https://github.com/openai/glide-text2im释放代码和权重。

In this work, we present a novel framework built to simplify 3D asset generation for amateur users. To enable interactive generation, our method supports a variety of input modalities that can be easily provided by a human, including images, text, partially observed shapes and combinations of these, further allowing to adjust the strength of each input. At the core of our approach is an encoder-decoder, compressing 3D shapes into a compact latent representation, upon which a diffusion model is learned. To enable a variety of multi-modal inputs, we employ task-specific encoders with dropout followed by a cross-attention mechanism. Due to its flexibility, our model naturally supports a variety of tasks, outperforming prior works on shape completion, image-based 3D reconstruction, and text-to-3D. Most interestingly, our model can combine all these tasks into one swiss-army-knife tool, enabling the user to perform shape generation using incomplete shapes, images, and textual descriptions at the same time, providing the relative weights for each input and facilitating interactivity. Despite our approach being shape-only, we further show an efficient method to texture the generated shape using large-scale text-to-image models.

While generative models produce high-quality images of concepts learned from a large-scale database, a user often wishes to synthesize instantiations of their own concepts (for example, their family, pets, or items). Can we teach a model to quickly acquire a new concept, given a few examples? Furthermore, can we compose multiple new concepts together? We propose Custom Diffusion, an efficient method for augmenting existing text-to-image models. We find that only optimizing a few parameters in the text-to-image conditioning mechanism is sufficiently powerful to represent new concepts while enabling fast tuning (~6 minutes). Additionally, we can jointly train for multiple concepts or combine multiple fine-tuned models into one via closed-form constrained optimization. Our fine-tuned model generates variations of multiple, new concepts and seamlessly composes them with existing concepts in novel settings. Our method outperforms several baselines and concurrent works, regarding both qualitative and quantitative evaluations, while being memory and computationally efficient.