文本对图像模型提供了前所未有的自由,可以通过自然语言指导创作。然而,尚不清楚如何行使这种自由以生成特定独特概念,修改其外观或以新角色和新颖场景构成它们的图像。换句话说,我们问:我们如何使用语言指导的模型将猫变成绘画,或者想象基于我们喜欢的玩具的新产品?在这里,我们提出了一种简单的方法,可以允许这种创造性自由。我们仅使用3-5个用户提供的概念(例如对象或样式)的图像,我们学会通过在冷冻文本到图像模型的嵌入空间中通过新的“单词”表示它。这些“单词”可以组成自然语言句子,以直观的方式指导个性化的创作。值得注意的是,我们发现有证据表明单词嵌入足以捕获独特而多样的概念。我们将我们的方法比较了各种基线,并证明它可以更忠实地描绘出一系列应用程序和任务的概念。我们的代码,数据和新单词将在以下网址提供:https://textual-inversion.github.io
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大型文本对图像模型在AI的演变中取得了显着的飞跃,从而使图像从给定的文本提示中实现了高质量和多样化的图像合成。但是,这些模型缺乏在给定的参考集中模仿受试者的外观,并在不同情况下合成它们的新颖性。在这项工作中,我们提出了一种新的方法,用于“个性化”文本图像扩散模型(将它们专门针对用户的需求)。仅作为一个主题的几张图像给出,我们将验证的文本对图像模型(图像,尽管我们的方法不限于特定模型),以便它学会了将唯一标识符与该特定主题结合。一旦将受试者嵌入模型的输出域中,就可以使用唯一标识符来合成主题的完全新颖的光真逼真的图像在不同场景中的上下文化。通过利用具有新的自动构基特异性的先前保存损失的语义先验嵌入到模型中,我们的技术可以在参考图像中未出现的不同场景,姿势,视图和照明条件中合成主题。我们将技术应用于几个以前无用的任务,包括主题重新定义,文本指导的视图合成,外观修改和艺术渲染(所有这些都保留了主题的关键特征)。项目页面:https://dreambooth.github.io/
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可以训练生成模型,以从特定域中生成图像,仅由文本提示引导,而不看到任何图像?换句话说:可以将图像生成器“盲目地训练”吗?利用大规模对比语言图像预训练(CLIP)模型的语义力量,我们提出了一种文本驱动方法,允许将生成模型转移到新域,而无需收集单个图像。我们展示通过自然语言提示和几分钟的培训,我们的方法可以通过各种风格和形状的多种域调整发电机。值得注意的是,许多这些修改难以与现有方法达到困难或完全不可能。我们在广泛的域中进行了广泛的实验和比较。这些证明了我们方法的有效性,并表明我们的移动模型保持了对下游任务吸引的生成模型的潜在空间属性。
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Recent works on diffusion models have demonstrated a strong capability for conditioning image generation, e.g., text-guided image synthesis. Such success inspires many efforts trying to use large-scale pre-trained diffusion models for tackling a challenging problem--real image editing. Works conducted in this area learn a unique textual token corresponding to several images containing the same object. However, under many circumstances, only one image is available, such as the painting of the Girl with a Pearl Earring. Using existing works on fine-tuning the pre-trained diffusion models with a single image causes severe overfitting issues. The information leakage from the pre-trained diffusion models makes editing can not keep the same content as the given image while creating new features depicted by the language guidance. This work aims to address the problem of single-image editing. We propose a novel model-based guidance built upon the classifier-free guidance so that the knowledge from the model trained on a single image can be distilled into the pre-trained diffusion model, enabling content creation even with one given image. Additionally, we propose a patch-based fine-tuning that can effectively help the model generate images of arbitrary resolution. We provide extensive experiments to validate the design choices of our approach and show promising editing capabilities, including changing style, content addition, and object manipulation. The code is available for research purposes at https://github.com/zhang-zx/SINE.git .
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While generative models produce high-quality images of concepts learned from a large-scale database, a user often wishes to synthesize instantiations of their own concepts (for example, their family, pets, or items). Can we teach a model to quickly acquire a new concept, given a few examples? Furthermore, can we compose multiple new concepts together? We propose Custom Diffusion, an efficient method for augmenting existing text-to-image models. We find that only optimizing a few parameters in the text-to-image conditioning mechanism is sufficiently powerful to represent new concepts while enabling fast tuning (~6 minutes). Additionally, we can jointly train for multiple concepts or combine multiple fine-tuned models into one via closed-form constrained optimization. Our fine-tuned model generates variations of multiple, new concepts and seamlessly composes them with existing concepts in novel settings. Our method outperforms several baselines and concurrent works, regarding both qualitative and quantitative evaluations, while being memory and computationally efficient.
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Large-scale diffusion-based generative models have led to breakthroughs in text-conditioned high-resolution image synthesis. Starting from random noise, such text-to-image diffusion models gradually synthesize images in an iterative fashion while conditioning on text prompts. We find that their synthesis behavior qualitatively changes throughout this process: Early in sampling, generation strongly relies on the text prompt to generate text-aligned content, while later, the text conditioning is almost entirely ignored. This suggests that sharing model parameters throughout the entire generation process may not be ideal. Therefore, in contrast to existing works, we propose to train an ensemble of text-to-image diffusion models specialized for different synthesis stages. To maintain training efficiency, we initially train a single model, which is then split into specialized models that are trained for the specific stages of the iterative generation process. Our ensemble of diffusion models, called eDiff-I, results in improved text alignment while maintaining the same inference computation cost and preserving high visual quality, outperforming previous large-scale text-to-image diffusion models on the standard benchmark. In addition, we train our model to exploit a variety of embeddings for conditioning, including the T5 text, CLIP text, and CLIP image embeddings. We show that these different embeddings lead to different behaviors. Notably, the CLIP image embedding allows an intuitive way of transferring the style of a reference image to the target text-to-image output. Lastly, we show a technique that enables eDiff-I's "paint-with-words" capability. A user can select the word in the input text and paint it in a canvas to control the output, which is very handy for crafting the desired image in mind. The project page is available at https://deepimagination.cc/eDiff-I/
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在Web规模数据上预测的大型视觉和语言模型提供了对许多V&L问题无价的表示。但是,目前尚不清楚如何将它们用于以非结构化语言为特定于用户特定的视觉概念。这个问题来自多个域,从个性化图像检索到与智能设备的个性化交互。我们介绍了一个新的学习设置,称为个性化视觉和语言(PERVL),并使用两个新的基准数据集来检索和细分用户特定的“个性化”概念“野外”。在PERVL中,应该独立于下游任务(2)允许经过审慎的模型以免费语言来推论它们,并且(3)不需要个性化的负面示例。我们提出了一个用于解决PERVL的体系结构,该体系结构通过扩展了一个预审计模型的输入词汇,并用新单词嵌入新的个性化概念。然后,模型可以通过简单地在句子中使用它们来推理它们。我们证明我们的方法从几个示例中学习了个性化的视觉概念,并且可以使用丰富的文本查询有效地将它们应用于图像检索和语义细分中。
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随着信息中的各种方式存在于现实世界中的各种方式,多式联信息之间的有效互动和融合在计算机视觉和深度学习研究中的多模式数据的创造和感知中起着关键作用。通过卓越的功率,在多式联运信息中建模互动,多式联运图像合成和编辑近年来已成为一个热门研究主题。与传统的视觉指导不同,提供明确的线索,多式联路指南在图像合成和编辑方面提供直观和灵活的手段。另一方面,该领域也面临着具有固有的模态差距的特征的几个挑战,高分辨率图像的合成,忠实的评估度量等。在本调查中,我们全面地阐述了最近多式联运图像综合的进展根据数据模型和模型架构编辑和制定分类。我们从图像合成和编辑中的不同类型的引导方式开始介绍。然后,我们描述了多模式图像综合和编辑方法,其具有详细的框架,包括生成的对抗网络(GAN),GaN反转,变压器和其他方法,例如NERF和扩散模型。其次是在多模式图像合成和编辑中广泛采用的基准数据集和相应的评估度量的综合描述,以及分析各个优点和限制的不同合成方法的详细比较。最后,我们为目前的研究挑战和未来的研究方向提供了深入了解。与本调查相关的项目可在HTTPS://github.com/fnzhan/mise上获得
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最近已被证明扩散模型产生高质量的合成图像,尤其是与指导技术配对,以促进忠诚的多样性。我们探索文本条件图像综合问题的扩散模型,并比较了两种不同的指导策略:剪辑指导和自由分类指导。我们发现后者是人类评估者的优选,用于光敏和标题相似度,并且通常产生光素质拟种样品。使用自由分类指导的35亿参数文本条件扩散模型的样本由人类评估者对来自Dall-E的人的人们青睐,即使后者使用昂贵的剪辑重新划分。此外,我们发现我们的模型可以进行微调,以执行图像修复,从而实现强大的文本驱动的图像编辑。我们在过滤的数据集中培训较小的模型,并在https://github.com/openai/glide-text2im释放代码和权重。
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利用深度学习的最新进展,文本到图像生成模型目前具有吸引公众关注的优点。其中两个模型Dall-E 2和Imagen已经证明,可以从图像的简单文本描述中生成高度逼真的图像。基于一种称为扩散模型的新型图像生成方法,文本对图像模型可以生产许多不同类型的高分辨率图像,其中人类想象力是唯一的极限。但是,这些模型需要大量的计算资源来训练,并处理从互联网收集的大量数据集。此外,代码库和模型均未发布。因此,它可以防止AI社区尝试这些尖端模型,从而使其结果复制变得复杂,即使不是不可能。在本文中,我们的目标是首先回顾这些模型使用的不同方法和技术,然后提出我们自己的文本模型模型实施。高度基于DALL-E 2,我们引入了一些轻微的修改,以应对所引起的高计算成本。因此,我们有机会进行实验,以了解这些模型的能力,尤其是在低资源制度中。特别是,我们提供了比Dall-e 2的作者(包括消融研究)更深入的分析。此外,扩散模型使用所谓的指导方法来帮助生成过程。我们引入了一种新的指导方法,该方法可以与其他指导方法一起使用,以提高图像质量。最后,我们的模型产生的图像质量相当好,而不必维持最先进的文本对图像模型的重大培训成本。
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我们介绍了自回归文本到图像(Parti)模型的途径,该模型生成高保真的影像图像并支持涉及复杂组成和世界知识的内容丰富的合成。 Parti将文本对图像生成视为类似于机器翻译的序列到序列建模问题,图像令牌的序列是目标输出,而不是其他语言的文本令牌。这种策略自然可以利用大型语言模型的先前工作,通过扩展数据和模型尺寸,能力和性能的持续进展。我们的方法很简单:首先,Parti使用基于变压器的图像令牌VIT-VQGAN将图像编码为离散令牌的序列。其次,我们通过将编码器二次变压器模型缩放到20B参数来实现一致的质量改进,其新的最新零弹药FID得分为7.23,而MS-Coco的FIDED得分为3.22。我们对本地化叙述以及党的详细分析(P2),这是1600多个英语提示的新的整体基准,证明了Parti在各种类别和难度方面的有效性。我们还探索并突出了我们的模型的局限性,以定义和体现关注重点领域以进一步改进。有关高分辨率图像,请参见https://parti.research.google/。
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我们提出了快速的文本2stylegan,这是一种自然语言界面,可适应预先训练的甘体,以实现文本引导的人脸合成。利用对比性语言图像预训练(剪辑)的最新进展,在培训过程中不需要文本数据。Fast Text2Stylegan被配制为条件变异自动编码器(CVAE),可在测试时为生成的图像提供额外的控制和多样性。我们的模型在遇到新的文本提示时不需要重新训练或微调剂或剪辑。与先前的工作相反,我们不依赖于测试时间的优化,这使我们的方法数量级比先前的工作快。从经验上讲,在FFHQ数据集上,我们的方法提供了与先前的工作相比,自然语言描述中具有不同详细程度的自然语言描述中的图像。
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We introduce M-VADER: a diffusion model (DM) for image generation where the output can be specified using arbitrary combinations of images and text. We show how M-VADER enables the generation of images specified using combinations of image and text, and combinations of multiple images. Previously, a number of successful DM image generation algorithms have been introduced that make it possible to specify the output image using a text prompt. Inspired by the success of those models, and led by the notion that language was already developed to describe the elements of visual contexts that humans find most important, we introduce an embedding model closely related to a vision-language model. Specifically, we introduce the embedding model S-MAGMA: a 13 billion parameter multimodal decoder combining components from an autoregressive vision-language model MAGMA and biases finetuned for semantic search.
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自然语言为图像编辑提供高度直观的界面。在本文中,我们基于自然语言描述与ROI掩模一起介绍用于在通用自然图像中执行局部(基于区域的)编辑的第一解决方案。我们通过利用并结合预先训练的语言图像模型(CLIP)来实现我们的目标,以使编辑朝向用户提供的文本提示,具有去噪扩散概率模型(DDPM)来产生自然的结果。为了使编辑区域与图像的不变部分无缝熔化,我们在噪声水平的进展下使用本地文本引导的扩散潜伏在空间上混合输入图像的声明版本。此外,我们表明向扩散过程增加增强,减轻了对抗性结果。我们与定性和定量的几个基线和相关方法进行比较,并表明我们的方法在整体现实主义方面优于这些解决方案,保留背景和匹配文本的能力。最后,我们显示了多个文本驱动的编辑应用程序,包括将新对象添加到图像,删除/替换/更改现有对象,背景替换和图像外推。
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最近,大规模文本驱动的合成模型由于其出色的产生高度多样化的图像而引起了很多关注,这些图像遵循给定的文本提示。这种基于文本的综合方法特别有吸引力,这些方法对人类用来口头描述其意图。因此,将文本驱动的图像合成扩展到文本驱动的图像编辑是很自然的。编辑对于这些生成模型来说是具有挑战性的,因为编辑技术的先天属性是保留大多数原始图像,而在基于文本的模型中,即使对文本提示的小修改也通常会导致完全不同的结果。最先进的方法可以通过要求用户提供空间掩码来本地化编辑,从而忽略蒙版区域内的原始结构和内容,从而减轻这种方法。在本文中,我们追求一个直观的及时提示编辑框架,其中编辑仅由文本控制。为此,我们深入分析了一个文本条件模型,并观察到跨注意层是控制图像的空间布局与提示中每个单词之间关系的关键。通过此观察,我们提出了几种应用程序,它们仅通过编辑文本提示来监视图像综合。这包括通过替换单词,通过添加规范来替换单词编辑的本地化编辑,甚至精心控制单词在图像中反映的程度。我们介绍了各种图像和提示的结果,证明了对编辑提示的高质量综合和忠诚度。
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从开放式域文本提示中生成和编辑图像是迄今为止需要昂贵且经过特殊训练的型号的一项挑战性的任务。我们为这两个任务展示了一种新颖的方法,该方法能够通过使用多模式编码器来指导图像世代,从而从具有显着语义复杂性的文本提示中产生高视觉质量的图像。我们在各种任务上说明了如何使用夹[37]引导VQGAN [11]产生的视觉质量输出比先前的较不灵活的方法,例如DALL-E [38],Glide [33]和Open-Edit [24],尽管没有接受培训的任务。我们的代码在公共存储库中可用。
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可控图像合成模型允许根据文本指令或来自示例图像的指导创建不同的图像。最近,已经显示出去噪扩散概率模型比现有方法产生更现实的图像,并且已在无条件和类条件设置中成功展示。我们探索细粒度,连续控制该模型类,并引入了一种新颖的统一框架,用于语义扩散指导,允许语言或图像指导,或两者。使用图像文本或图像匹配分数的梯度将指导注入预训练的无条件扩散模型中。我们探讨基于剪辑的文本指导,以及以统一形式的基于内容和类型的图像指导。我们的文本引导综合方法可以应用于没有相关文本注释的数据集。我们对FFHQ和LSUN数据集进行实验,并显示出细粒度的文本引导图像合成的结果,与样式或内容示例图像相关的图像的合成,以及具有文本和图像引导的示例。
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2D-to-3D reconstruction is an ill-posed problem, yet humans are good at solving this problem due to their prior knowledge of the 3D world developed over years. Driven by this observation, we propose NeRDi, a single-view NeRF synthesis framework with general image priors from 2D diffusion models. Formulating single-view reconstruction as an image-conditioned 3D generation problem, we optimize the NeRF representations by minimizing a diffusion loss on its arbitrary view renderings with a pretrained image diffusion model under the input-view constraint. We leverage off-the-shelf vision-language models and introduce a two-section language guidance as conditioning inputs to the diffusion model. This is essentially helpful for improving multiview content coherence as it narrows down the general image prior conditioned on the semantic and visual features of the single-view input image. Additionally, we introduce a geometric loss based on estimated depth maps to regularize the underlying 3D geometry of the NeRF. Experimental results on the DTU MVS dataset show that our method can synthesize novel views with higher quality even compared to existing methods trained on this dataset. We also demonstrate our generalizability in zero-shot NeRF synthesis for in-the-wild images.
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We extend Textual Inversion to learn pseudo-words that represent a concept at different resolutions. This allows us to generate images that use the concept with different levels of detail and also to manipulate different resolutions using language. Once learned, the user can generate images at different levels of agreement to the original concept; "A photo of $S^*(0)$" produces the exact object while the prompt "A photo of $S^*(0.8)$" only matches the rough outlines and colors. Our framework allows us to generate images that use different resolutions of an image (e.g. details, textures, styles) as separate pseudo-words that can be composed in various ways. We open-soure our code in the following URL: https://github.com/giannisdaras/multires_textual_inversion
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Recent 3D generative models have achieved remarkable performance in synthesizing high resolution photorealistic images with view consistency and detailed 3D shapes, but training them for diverse domains is challenging since it requires massive training images and their camera distribution information. Text-guided domain adaptation methods have shown impressive performance on converting the 2D generative model on one domain into the models on other domains with different styles by leveraging the CLIP (Contrastive Language-Image Pre-training), rather than collecting massive datasets for those domains. However, one drawback of them is that the sample diversity in the original generative model is not well-preserved in the domain-adapted generative models due to the deterministic nature of the CLIP text encoder. Text-guided domain adaptation will be even more challenging for 3D generative models not only because of catastrophic diversity loss, but also because of inferior text-image correspondence and poor image quality. Here we propose DATID-3D, a domain adaptation method tailored for 3D generative models using text-to-image diffusion models that can synthesize diverse images per text prompt without collecting additional images and camera information for the target domain. Unlike 3D extensions of prior text-guided domain adaptation methods, our novel pipeline was able to fine-tune the state-of-the-art 3D generator of the source domain to synthesize high resolution, multi-view consistent images in text-guided targeted domains without additional data, outperforming the existing text-guided domain adaptation methods in diversity and text-image correspondence. Furthermore, we propose and demonstrate diverse 3D image manipulations such as one-shot instance-selected adaptation and single-view manipulated 3D reconstruction to fully enjoy diversity in text.
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