In this paper, we propose a diffusion-based face swapping framework for the first time, called DiffFace, composed of training ID conditional DDPM, sampling with facial guidance, and a target-preserving blending. In specific, in the training process, the ID conditional DDPM is trained to generate face images with the desired identity. In the sampling process, we use the off-the-shelf facial expert models to make the model transfer source identity while preserving target attributes faithfully. During this process, to preserve the background of the target image and obtain the desired face swapping result, we additionally propose a target-preserving blending strategy. It helps our model to keep the attributes of the target face from noise while transferring the source facial identity. In addition, without any re-training, our model can flexibly apply additional facial guidance and adaptively control the ID-attributes trade-off to achieve the desired results. To the best of our knowledge, this is the first approach that applies the diffusion model in face swapping task. Compared with previous GAN-based approaches, by taking advantage of the diffusion model for the face swapping task, DiffFace achieves better benefits such as training stability, high fidelity, diversity of the samples, and controllability. Extensive experiments show that our DiffFace is comparable or superior to the state-of-the-art methods on several standard face swapping benchmarks.
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最近,GaN反演方法与对比语言 - 图像预先绘制(CLIP)相结合,可以通过文本提示引导零拍摄图像操作。然而,由于GaN反转能力有限,它们对不同实物的不同实物的应用仍然困难。具体地,这些方法通常在与训练数据相比,改变对象标识或产生不需要的图像伪影的比较与新颖姿势,视图和高度可变内容重建具有新颖姿势,视图和高度可变内容的困难。为了减轻这些问题并实现真实图像的忠实操纵,我们提出了一种新的方法,Dumbused Clip,其使用扩散模型执行文本驱动的图像操纵。基于近期扩散模型的完整反转能力和高质量的图像生成功率,即使在看不见的域之间也成功地执行零拍摄图像操作。此外,我们提出了一种新颖的噪声组合方法,允许简单的多属性操作。与现有基线相比,广泛的实验和人类评估确认了我们的方法的稳健和卓越的操纵性能。
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Inspired by the impressive performance of recent face image editing methods, several studies have been naturally proposed to extend these methods to the face video editing task. One of the main challenges here is temporal consistency among edited frames, which is still unresolved. To this end, we propose a novel face video editing framework based on diffusion autoencoders that can successfully extract the decomposed features - for the first time as a face video editing model - of identity and motion from a given video. This modeling allows us to edit the video by simply manipulating the temporally invariant feature to the desired direction for the consistency. Another unique strength of our model is that, since our model is based on diffusion models, it can satisfy both reconstruction and edit capabilities at the same time, and is robust to corner cases in wild face videos (e.g. occluded faces) unlike the existing GAN-based methods.
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Diffusion-based generative models have achieved remarkable success in image generation. Their guidance formulation allows an external model to plug-and-play control the generation process for various tasks without fine-tuning the diffusion model. However, the direct use of publicly available off-the-shelf models for guidance fails due to their poor performance on noisy inputs. For that, the existing practice is to fine-tune the guidance models with labeled data corrupted with noises. In this paper, we argue that this practice has limitations in two aspects: (1) performing on inputs with extremely various noises is too hard for a single model; (2) collecting labeled datasets hinders scaling up for various tasks. To tackle the limitations, we propose a novel strategy that leverages multiple experts where each expert is specialized in a particular noise range and guides the reverse process at its corresponding timesteps. However, as it is infeasible to manage multiple networks and utilize labeled data, we present a practical guidance framework termed Practical Plug-And-Play (PPAP), which leverages parameter-efficient fine-tuning and data-free knowledge transfer. We exhaustively conduct ImageNet class conditional generation experiments to show that our method can successfully guide diffusion with small trainable parameters and no labeled data. Finally, we show that image classifiers, depth estimators, and semantic segmentation models can guide publicly available GLIDE through our framework in a plug-and-play manner.
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可控图像合成模型允许根据文本指令或来自示例图像的指导创建不同的图像。最近,已经显示出去噪扩散概率模型比现有方法产生更现实的图像,并且已在无条件和类条件设置中成功展示。我们探索细粒度,连续控制该模型类,并引入了一种新颖的统一框架,用于语义扩散指导,允许语言或图像指导,或两者。使用图像文本或图像匹配分数的梯度将指导注入预训练的无条件扩散模型中。我们探讨基于剪辑的文本指导,以及以统一形式的基于内容和类型的图像指导。我们的文本引导综合方法可以应用于没有相关文本注释的数据集。我们对FFHQ和LSUN数据集进行实验,并显示出细粒度的文本引导图像合成的结果,与样式或内容示例图像相关的图像的合成,以及具有文本和图像引导的示例。
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在本文中,我们解决了神经面部重演的问题,鉴于一对源和目标面部图像,我们需要通过将目标的姿势(定义为头部姿势及其面部表情定义)通过同时保留源的身份特征(例如面部形状,发型等),即使在源头和目标面属于不同身份的挑战性情况下也是如此。在此过程中,我们解决了最先进作品的一些局限在推理期间标记的数据以及c)它们不保留大型头部姿势变化中的身份。更具体地说,我们提出了一个框架,该框架使用未配对的随机生成的面部图像学会通过合并最近引入的样式空间$ \ Mathcal $ \ Mathcal {S} $ of Stylegan2的姿势,以将面部的身份特征从其姿势中解脱出来表现出显着的分解特性。通过利用这一点,我们学会使用3D模型的监督成功地混合了一对源和目标样式代码。随后用于重新制定的最终潜在代码由仅与源的面部姿势相对应的潜在单位和仅与源身份相对应的单位组成,从而显着改善了与最近的状态性能相比的重新制定性能。艺术方法。与艺术的状态相比,我们定量和定性地表明,即使在极端的姿势变化下,提出的方法也会产生更高的质量结果。最后,我们通过首先将它们嵌入预告片发电机的潜在空间来报告实际图像。我们在:https://github.com/stelabou/stylemask上公开提供代码和预估计的模型
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鉴于其广泛的应用,已经对人面部交换的任务进行了许多尝试。尽管现有的方法主要依赖于乏味的网络和损失设计,但它们仍然在源和目标面之间的信息平衡中挣扎,并倾向于产生可见的人工制品。在这项工作中,我们引入了一个名为StylesWap的简洁有效的框架。我们的核心想法是利用基于样式的生成器来增强高保真性和稳健的面部交换,因此可以采用发电机的优势来优化身份相似性。我们仅通过最小的修改来确定,StyleGAN2体系结构可以成功地处理来自源和目标的所需信息。此外,受到TORGB层的启发,进一步设计了交换驱动的面具分支以改善信息的融合。此外,可以采用stylegan倒置的优势。特别是,提出了交换引导的ID反转策略来优化身份相似性。广泛的实验验证了我们的框架会产生高质量的面部交换结果,从而超过了最先进的方法,既有定性和定量。
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Conditional diffusion probabilistic models can model the distribution of natural images and can generate diverse and realistic samples based on given conditions. However, oftentimes their results can be unrealistic with observable color shifts and textures. We believe that this issue results from the divergence between the probabilistic distribution learned by the model and the distribution of natural images. The delicate conditions gradually enlarge the divergence during each sampling timestep. To address this issue, we introduce a new method that brings the predicted samples to the training data manifold using a pretrained unconditional diffusion model. The unconditional model acts as a regularizer and reduces the divergence introduced by the conditional model at each sampling step. We perform comprehensive experiments to demonstrate the effectiveness of our approach on super-resolution, colorization, turbulence removal, and image-deraining tasks. The improvements obtained by our method suggest that the priors can be incorporated as a general plugin for improving conditional diffusion models.
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与生成的对抗网(GAN)相比,降级扩散概率模型(DDPM)在各种图像生成任务中取得了显着成功。关于语义图像综合的最新工作主要遵循\ emph {de exto}基于gan的方法,这可能导致生成图像的质量或多样性不令人满意。在本文中,我们提出了一个基于DDPM的新型框架,用于语义图像合成。与先前的条件扩散模型不同,将语义布局和嘈杂的图像作为输入为U-NET结构,该结构可能无法完全利用输入语义掩码中的信息,我们的框架处理语义布局和嘈杂的图像不同。它将噪声图像馈送到U-NET结构的编码器时,而语义布局通过多层空间自适应归一化操作符将语义布局馈送到解码器。为了进一步提高语义图像合成中的发电质量和语义解释性,我们介绍了无分类器的指导采样策略,该策略承认采样过程的无条件模型的得分。在三个基准数据集上进行的广泛实验证明了我们提出的方法的有效性,从而在忠诚度(FID)和多样性〜(LPIPS)方面实现了最先进的性能。
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Object compositing based on 2D images is a challenging problem since it typically involves multiple processing stages such as color harmonization, geometry correction and shadow generation to generate realistic results. Furthermore, annotating training data pairs for compositing requires substantial manual effort from professionals, and is hardly scalable. Thus, with the recent advances in generative models, in this work, we propose a self-supervised framework for object compositing by leveraging the power of conditional diffusion models. Our framework can hollistically address the object compositing task in a unified model, transforming the viewpoint, geometry, color and shadow of the generated object while requiring no manual labeling. To preserve the input object's characteristics, we introduce a content adaptor that helps to maintain categorical semantics and object appearance. A data augmentation method is further adopted to improve the fidelity of the generator. Our method outperforms relevant baselines in both realism and faithfulness of the synthesized result images in a user study on various real-world images.
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In recent years, generative models have undergone significant advancement due to the success of diffusion models. The success of these models is often attributed to their use of guidance techniques, such as classifier and classifier-free methods, which provides effective mechanisms to trade-off between fidelity and diversity. However, these methods are not capable of guiding a generated image to be aware of its geometric configuration, e.g., depth, which hinders the application of diffusion models to areas that require a certain level of depth awareness. To address this limitation, we propose a novel guidance approach for diffusion models that uses estimated depth information derived from the rich intermediate representations of diffusion models. To do this, we first present a label-efficient depth estimation framework using the internal representations of diffusion models. At the sampling phase, we utilize two guidance techniques to self-condition the generated image using the estimated depth map, the first of which uses pseudo-labeling, and the subsequent one uses a depth-domain diffusion prior. Experiments and extensive ablation studies demonstrate the effectiveness of our method in guiding the diffusion models toward geometrically plausible image generation. Project page is available at https://ku-cvlab.github.io/DAG/.
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近年来,面部交换模型在发电质量方面取得了进展,并引起了他们在隐私保护和娱乐中的应用。然而,他们的复杂架构和损失函数通常需要仔细调整成功培训。在本文中,我们提出了一种名为“顺利交换”的新型脸部交换模型,专注于导出身份嵌入的平滑度,而不是采用复杂的手工设计。我们假设面部交换难度的要点是不稳定的梯度,并且可以通过平滑的身份嵌入器来解决。顺利交换采用使用监督对比学习的嵌入式培训,在那里我们发现其改进的平滑度允许更快,稳定的培训,即使是基于简单的基于U-Net的发电机和三个基本损耗功能。在野外的面部交换基准(FFHQ,面部素材++)和面部图像的广泛实验表明,在身份变化方面,我们的模型也定量和质量上的可比性甚至优于现有的方法。
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Recent works on diffusion models have demonstrated a strong capability for conditioning image generation, e.g., text-guided image synthesis. Such success inspires many efforts trying to use large-scale pre-trained diffusion models for tackling a challenging problem--real image editing. Works conducted in this area learn a unique textual token corresponding to several images containing the same object. However, under many circumstances, only one image is available, such as the painting of the Girl with a Pearl Earring. Using existing works on fine-tuning the pre-trained diffusion models with a single image causes severe overfitting issues. The information leakage from the pre-trained diffusion models makes editing can not keep the same content as the given image while creating new features depicted by the language guidance. This work aims to address the problem of single-image editing. We propose a novel model-based guidance built upon the classifier-free guidance so that the knowledge from the model trained on a single image can be distilled into the pre-trained diffusion model, enabling content creation even with one given image. Additionally, we propose a patch-based fine-tuning that can effectively help the model generate images of arbitrary resolution. We provide extensive experiments to validate the design choices of our approach and show promising editing capabilities, including changing style, content addition, and object manipulation. The code is available for research purposes at https://github.com/zhang-zx/SINE.git .
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Generative models, particularly GANs, have been utilized for image editing. Although GAN-based methods perform well on generating reasonable contents aligned with the user's intentions, they struggle to strictly preserve the contents outside the editing region. To address this issue, we use diffusion models instead of GANs and propose a novel image-editing method, based on pixel-wise guidance. Specifically, we first train pixel-classifiers with few annotated data and then estimate the semantic segmentation map of a target image. Users then manipulate the map to instruct how the image is to be edited. The diffusion model generates an edited image via guidance by pixel-wise classifiers, such that the resultant image aligns with the manipulated map. As the guidance is conducted pixel-wise, the proposed method can create reasonable contents in the editing region while preserving the contents outside this region. The experimental results validate the advantages of the proposed method both quantitatively and qualitatively.
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数字艺术合成在多媒体社区中受到越来越多的关注,因为有效地与公众参与了艺术。当前的数字艺术合成方法通常使用单模式输入作为指导,从而限制了模型的表现力和生成结果的多样性。为了解决这个问题,我们提出了多模式引导的艺术品扩散(MGAD)模型,该模型是一种基于扩散的数字艺术品生成方法,它利用多模式提示作为控制无分类器扩散模型的指导。此外,对比度语言图像预处理(剪辑)模型用于统一文本和图像模式。关于生成的数字艺术绘画质量和数量的广泛实验结果证实了扩散模型和多模式指导的组合有效性。代码可从https://github.com/haha-lisa/mgad-multimodal-guided-artwork-diffusion获得。
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Generic image inpainting aims to complete a corrupted image by borrowing surrounding information, which barely generates novel content. By contrast, multi-modal inpainting provides more flexible and useful controls on the inpainted content, \eg, a text prompt can be used to describe an object with richer attributes, and a mask can be used to constrain the shape of the inpainted object rather than being only considered as a missing area. We propose a new diffusion-based model named SmartBrush for completing a missing region with an object using both text and shape-guidance. While previous work such as DALLE-2 and Stable Diffusion can do text-guided inapinting they do not support shape guidance and tend to modify background texture surrounding the generated object. Our model incorporates both text and shape guidance with precision control. To preserve the background better, we propose a novel training and sampling strategy by augmenting the diffusion U-net with object-mask prediction. Lastly, we introduce a multi-task training strategy by jointly training inpainting with text-to-image generation to leverage more training data. We conduct extensive experiments showing that our model outperforms all baselines in terms of visual quality, mask controllability, and background preservation.
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Generative models have been widely studied in computer vision. Recently, diffusion models have drawn substantial attention due to the high quality of their generated images. A key desired property of image generative models is the ability to disentangle different attributes, which should enable modification towards a style without changing the semantic content, and the modification parameters should generalize to different images. Previous studies have found that generative adversarial networks (GANs) are inherently endowed with such disentanglement capability, so they can perform disentangled image editing without re-training or fine-tuning the network. In this work, we explore whether diffusion models are also inherently equipped with such a capability. Our finding is that for stable diffusion models, by partially changing the input text embedding from a neutral description (e.g., "a photo of person") to one with style (e.g., "a photo of person with smile") while fixing all the Gaussian random noises introduced during the denoising process, the generated images can be modified towards the target style without changing the semantic content. Based on this finding, we further propose a simple, light-weight image editing algorithm where the mixing weights of the two text embeddings are optimized for style matching and content preservation. This entire process only involves optimizing over around 50 parameters and does not fine-tune the diffusion model itself. Experiments show that the proposed method can modify a wide range of attributes, with the performance outperforming diffusion-model-based image-editing algorithms that require fine-tuning. The optimized weights generalize well to different images. Our code is publicly available at https://github.com/UCSB-NLP-Chang/DiffusionDisentanglement.
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Generating photos satisfying multiple constraints find broad utility in the content creation industry. A key hurdle to accomplishing this task is the need for paired data consisting of all modalities (i.e., constraints) and their corresponding output. Moreover, existing methods need retraining using paired data across all modalities to introduce a new condition. This paper proposes a solution to this problem based on denoising diffusion probabilistic models (DDPMs). Our motivation for choosing diffusion models over other generative models comes from the flexible internal structure of diffusion models. Since each sampling step in the DDPM follows a Gaussian distribution, we show that there exists a closed-form solution for generating an image given various constraints. Our method can unite multiple diffusion models trained on multiple sub-tasks and conquer the combined task through our proposed sampling strategy. We also introduce a novel reliability parameter that allows using different off-the-shelf diffusion models trained across various datasets during sampling time alone to guide it to the desired outcome satisfying multiple constraints. We perform experiments on various standard multimodal tasks to demonstrate the effectiveness of our approach. More details can be found in https://nithin-gk.github.io/projectpages/Multidiff/index.html
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Can a text-to-image diffusion model be used as a training objective for adapting a GAN generator to another domain? In this paper, we show that the classifier-free guidance can be leveraged as a critic and enable generators to distill knowledge from large-scale text-to-image diffusion models. Generators can be efficiently shifted into new domains indicated by text prompts without access to groundtruth samples from target domains. We demonstrate the effectiveness and controllability of our method through extensive experiments. Although not trained to minimize CLIP loss, our model achieves equally high CLIP scores and significantly lower FID than prior work on short prompts, and outperforms the baseline qualitatively and quantitatively on long and complicated prompts. To our best knowledge, the proposed method is the first attempt at incorporating large-scale pre-trained diffusion models and distillation sampling for text-driven image generator domain adaptation and gives a quality previously beyond possible. Moreover, we extend our work to 3D-aware style-based generators and DreamBooth guidance.
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我们提出了一种基于示例的图像翻译的新方法,称为匹配交织的扩散模型(MIDMS)。该任务的大多数现有方法都是基于GAN的匹配,然后代表了代代框架。但是,在此框架中,跨跨域的语义匹配难度引起的匹配误差,例如草图和照片,可以很容易地传播到生成步骤,从而导致结果退化。由于扩散模型的最新成功激发了克服GAN的缺点,我们结合了扩散模型以克服这些局限性。具体而言,我们制定了一个基于扩散的匹配和生成框架,该框架通过将中间扭曲馈入尖锐的过程并将其变形以生成翻译的图像,从而交织了潜在空间中的跨域匹配和扩散步骤。此外,为了提高扩散过程的可靠性,我们使用周期一致性设计了一种置信度的过程,以在翻译过程中仅考虑自信区域。实验结果表明,我们的MIDM比最新方法产生的图像更合理。
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